A simple proof of a theorem of Kronecker

Kronecker's theorem of the title is the following.

Theorem. — Let $\alpha\in\mathbf C$ be an algebraic integer all of whose conjugates have absolute value at most $1$. Then either $\alpha=0$, or $\alpha$ is a root of unity.

This theorem has several elementary proofs. In this post, I explain the simple proof proposed by Gebhart Greiter in his American Mathematical Monthly note, adding details so that it would (hopefully) be more accessible to students.

The possibility $\alpha=0$ is rather unentertaining, hence let us assume that $\alpha\neq0$.

Let us first analyse the hypothesis. The assumption that $\alpha$ is an algebraic integer means that there exists a monic polynomial $f\in\mathbf Z[T]$ such that $f(\alpha)=0$. I claim that $f$ can be assumed to be irreducible in $\mathbf Q[T]$; it is then the minimal polynomial of $\alpha$. This follows from Gauss's theorem, and let me give a proof for that. Let $g\in \mathbf Q[T]$ be a monic irreducible factor of $f$ and let $h\in\mathbf Q[T]$ such that $f=gh$. Chasing denominators and dividing out by their gcd, there are polynomials $g_1,h_1\in\mathbf Z[T]$ whose coefficients are mutually coprime, and natural integers $u,v$ such that $g=ug_1$ and $h=vh_1$. Then $(uv)f=g_1 h_1$. Since $f$ is monic, this implies that $uv$ is an integer. Let us prove that $uv=1$; otherwise, it has a prime factor $p$. Considering the relation $(uv)f=g_1h_1$ modulo $p$ gives $0=g_1 h_1 \pmod p$. Since their coefficients are mutually coprime, the polynomials $g_1$ and $h_1$ are nonzero modulo $p$, hence their product is nonzero. This is a contradiction.

So we have a monic polynomial $f\in\mathbf Z[T]$, irreducible in $\mathbf Q[T]$, such that $f(\alpha)=0$. That is to say, $f$ is the minimal polynomial of $\alpha$, so that the conjugates of $\alpha$ are the roots of $f$. Note that the roots of $f$ are pairwise distinct — otherwise, the $\gcd(f,f')$ would be a nontrivial factor of $f$. Moreover, $0$ is not a root of $f$, for otherwhise one could factor $f=T\cdot f_1$.

Let us now consider the companion matrix to $f$: writing $f=T^n+c_1T^{n-1}+\dots+c_n$, so that $n=\deg(f)$, this is the matrix \[ C_f = \begin{pmatrix} 0 & \dots & 0 & -c_n \\ 1 & \ddots & \vdots & -c_{n-1} \\ 0 & \ddots & & -c_{n-2} \\ 0 & \cdots & 1 & -c_1 \end{pmatrix}.\] If $e_1,\ldots,e_n$ are the canonical column vectors $e_1=(1,0,\dots,0)$, etc., then $C_f\cdot e_1=e_2$, \ldots, $C_f \cdot e_{n-1}=e_{n}$, and $C_f\cdot e_n = -c_{n} e_1-\dots -c_1 e_n$. Consequently, \[ f(C_f)\cdot e_1 = C_f^n\cdot e_1 +c_1 C_f^{n_1} \cdot e_1+ \dots +c_n e_1 = 0.\] Moreover, for $2\leq k\leq n$, one has $f(C_f)\cdot e_k=f(C_f)\cdot C_f^{k-1}\cdot e_1=C_f^{k-1}\cdot f(C_f)\cdot e_1=0$. Consequently, $f(C_f)=0$ and the complex eigenvalues of $f(C_f)$ are roots of $f$. Since $f$ has simple roots, $C_f$ is diagonalizable and their exists a matrix $P\in\mathrm{GL}(n,\mathbf C)$ and diagonal matrix $D$ such that $C_f=P\cdot D\cdot P^{-1}$, the diagonal entries of $D$ are roots of $f$. Since $0$ is not a root of $f$, the matrix $D$ is invertible, and $C_f$ is invertible as well. More precisely, one can infer from the definition of $C_f$ that $g(C_f)\cdot e_1\neq 0$ for any nonzero polynomial $g$ of degre $<n$, so that $f$ is the minimal polynomial of $C_f$. Consequently, all of its roots are actually eigenvalues of $C_f$, and appear in $D$; in particular, $\alpha$ is an eigenvalue of $C_f$.

For every $k\geq 1$, one has $C_f^k=P\cdot (D^k)\cdot P^{-1}$. Since all entries of $D$ have absolute value at most $1,$ the set of all $D^k$ is bounded in $\mathrm{M}_n(\mathbf C)$. Consequently, the set $\{C_f^k\,;\, k\in\mathbf Z\}$ is bounded in $\mathrm M_n(\mathbf C)$. On the other hand, this set consists in matrices in $\mathrm M_n(\mathbf Z)$. It follows that this set is finite.

There are integers $k$ and $\ell$ such that $k< \ell$ and $C_f^k=C_f^\ell$. Since $C_f$ is invertible, one has $C_f^{\ell-k}=1$. Since $\alpha$ is an eigenvalue of $C_f,$ this implies $\alpha^{\ell-k}=1$. We thus have proved that $\alpha$ is a root of unity.

Combinatorics of the nilpotent cone

$\global\def\Card{\operatorname{Card}}\global\def\GL{\mathrm{GL}}\global\def\im{\operatorname{im}}\gdef\KVar{\mathrm{KVar}}$

Let $n$ be an integer and $F$ be a field. Nilpotent matrices $N\in \mathrm M_n(F)$ are those matrices for which there exists an integer $p$ with $N^p=0$. Their characteristic polynomial is $\chi_N(T)=T^n$, and they satisfy $N^n=0$, which shows that the set $\mathscr N_n$ of nilpotent matrices is an algebraic variety. The equation $N^n=0$ is homogeneous of degree $n$, so that $\mathscr N_n$ is a cone.

The classification of nilpotent matrices is an intermediate step in the theory of Jordan decomposition: In an adequate basis, a nilpotent matrix can be written as a diagonal block matrix of “basic” nilpotent matrices, $p \times p$ matrices of the form \[ \begin{pmatrix} 0 & 0 & \dots & 0 & 0 \\ 1 & 0 & & & \vdots \\ 0 & 1 & \ddots & & 0 \\ \vdots & \ddots & \ddots & \ddots & 0 \\ 0 & & 0 & 1 & 0\end{pmatrix} \] whose minimal polynomial is $T^p$. The sum of the sizes of these blocks is $n$ and in this way, it is associated with any $n\times n$ nilpotent matrix a partition $\pi$ of~$n$. It is known that two nilpotent matrices are conjugate if and only if they are associated with the same partition. For any partition $\pi$ of~$n$, let us denote by $J_\pi$ the corresponding matrix whose sizes of blocks are arranged in increasing order, and $\mathscr N_\pi$ the set of nilpotent matrices that are associated with the partition $\pi$.

The theorem of Fine and Herstein (1958)

Having to teach “agrégation” classes made me learn about a classic combinatorial result: counting the number of nilpotent matrices when $F$ is a finite field.

Theorem (Fine, Herstein, 1958). — Let $F$ be a finite field with $q$ elements. The cardinality of $\mathscr N_n(F)$ is $q^{n^2-n}$. Equivalently, the probability that an $n\times n$ matrix with coefficients in $F$ be nilpotent is $q^{-n}$.

The initial proof of this results relies on the action of $\GL_n(F)$ on $\mathscr N_n(F)$: we recalled that the orbits correspond with the partitions of $n$, hence a decomposition \[ \Card(\mathscr N_n(F)) = \sum_{\pi} \Card(\mathscr N_\pi(F)). \] We know that $\mathscr N_\pi(F)$ is the orbit of the matrix $J_\pi$ under the action of $\GL_n(F)$. By the classic orbit-stabilizer formula, one thus has \[ \Card(\mathscr N_\pi(F)) = \frac{\Card(\GL_n(F))}{\Card(C_\pi(F))}, \] where $C_\pi(F)$ is the set of matrices $A\in\GL_n(F)$ such that $AJ_\pi=J_\pi A$. The precise description of $C_\pi(F)$ is delicate but their arguments go as follow.

They first replace the group $C_\pi(F)$ by the algebra $A_\pi(F)$ of all matrices $A\in\mathrm M_n(F)$ such that $AJ_\pi=J_\pi A$. For any integer, let $m_i$ be the multiplicity of an integer $i$ in the partition $\pi$, so that $n=\sum i m_i$. The block decomposition of $J_\pi$ corresponds with a decomposition of $F^n$ as a direct sum of invariant subspaces $V_i$, where $V_i$ has dimension $i m_i$. In fact, $V_1=\ker(J_\pi)$, $V_1\oplus V_2=\ker(J_\pi^2)$, etc. This shows that $A_\pi(F)$ is an algebra of block-triangular matrices. Moreover, the possible diagonal blocks can be shown to be isomorphic to $\mathrm M_{m_i}(F)$. In other words, we have a surjective morphism of algebras \[ A_\pi(F) \to \prod_i \mathrm M_{m_i}(F), \] whose kernel consists of nilpotent matrices. In particular, the proportion of invertible elements in $A_\pi(F)$ is equal to the proportion of invertible elements in the product $\prod_i \mathrm M_{m_i}(F)$.

Ultimately, Fine and Herstein obtain an explicit sum over the set of partitions of $n$ which they prove equals $q^{n^2-n}$, after an additional combinatorial argument.

Soon after, the theorem of Fine and Herstein was given easier proofs, starting from Gerstenhaber (1961) to Kaplansky (1990) and Leinster (2021).

A proof

The following proof is borrowed from Caldero and Peronnier (2022), Carnet de voyage en Algébrie. It can be seen as a simplification of the proofs of Gerstenhaber (1961) and Leinster (2021).

Let us start with the Fitting decomposition of an endomorphism $u\in \mathrm N_n(F)$: the least integer $p$ such that $\ker(u^p)=\ker(u^{p+1})$ coincides with the least integer $p$ such that $\im(u^p)=\im(u^{p+1})$, and one has $F^n=\ker(u^p)\oplus \im(u^p)$. The subspaces $N(u)=\ker(u^p)$ and $I(u)=\im(u^p)$ are invariant under $u$, and $u$ acts nilpotently on $\ker(u^p)$ and bijectively on $\im(u^p)$. In other words, we have associated with $u$ complementary subspaces $N(u)$ and $I(u)$, a nilpotent operator of $N(u)$ and an invertible operator on $I(u)$. This map is bijective.

For any integer $d$, let $\nu_d$ be the cardinality of nilpotent matrices in $\mathrm M_d(F)$, and $\gamma_d$ be the cardinality of invertible matrices in $\mathrm M_d(F)$. Let also $\mathscr D_d$ be the set of all pairs $(K,I)$, where $K$ and $I$ are complementary subspaces of dimensions $d$, $n-d$ of $F^n$. We thus obtain \[ n^2 = \sum_{(K,I)\in\mathscr D_d} \nu_d \cdot \gamma_{n-d}. \] We need to compute the cardinality of $\mathscr D_d$. In fact, given one pair $(K,I)\in\mathscr D_d$, all other are of the form $(g\cdot K,g\cdot I)$, for some $g\in\GL_n(F)$: the group $\GL_n(F)$ acts transitively on $\mathscr D_d$. The stabilizer of $(K,I)$ can be identified with $\GL_d(F)\times \GL_{n-d}(F)$. Consequently, \[ \Card(\mathscr D_d) = \frac{\Card(\GL_n(F))}{\Card(\GL_d(F)\Card(\GL_{n-d}(F))} = \frac{\gamma_n}{\gamma_d \gamma_{n-d}}. \] We thus obtain \[ q^{n^2} = \sum_{d=0}^n \frac{\gamma_n}{\gamma_d \gamma_{n-d}} \nu_d \gamma_{n-d} = \gamma_n \sum_{d=0}^n \frac{\nu_d}{\gamma_d}. \] By subtraction, we get \[ \frac{\nu_n}{\gamma_n} = \frac {q^{n^2}}{\gamma_n} - \frac{q^{(n-1)^2}}{\gamma_{n-1}},\] or \[ \nu_n = q^{n^2} - q^{(n-1)^2} \frac{\gamma_n}{\gamma_{n-1}}. \] It remains to compute $\gamma_n$: since an invertible matrix consists of a nonzero vector, a vector which does not belong to the line generated by the first one, etc., we have \[ \gamma_n = (q^n-1) (q^n-q)\dots (q^n-q^{n-1}). \] Then, \[ \gamma_n = (q^n-1) q^{n-1} (q^{n-1}-1)\dots (q^{n-1}-q^{n-2}) = (q^n-1)q^{n-1} \gamma_{n-1}. \] We thus obtain \[ \nu_n = q^{n^2} - q^{(n-1)^2} (q^n-1) q^{n-1} = q^{n^2} - q^{(n-1)^2} q^{2n-1} + q^{(n-1)^2} q^{n-1} = q^{n^2-n}, \] as claimed.

The proof of Leinster (2021)

Leinster defines a bijection from $\mathscr N_n(F)\times F^n$ to $\mathrm M_n(F)$. The definition is however not very canonical, because he assumes given, for any subspace $V$ of $F^n$, a basis of $V$.

Take a pair $(u,x)$, where $u\in\mathscr N_n(F)$ and $x\in F^n$ and consider the subspace $V_x=\langle x,u(x),\dots\rangle$, the smallest $u$-invariant subspace of $F^n$ which contains $x$. To describe $u$, we observe that we know its restriction to $V_x$, and we need to describe it on the chosen complementary subspace $V'_x$.

To that aim, we have to give ourselves an endomorphism $u'_x$ of $V'_x$ and a linear map $\phi_x\colon V'_x\to V_x$. Since we want $u$ to be nilpotent, it is necessary and sufficient to take $u'_x$ nilpotent.

Instead of considering $\phi_x\colon V'_x\to V_x$, we can consider the map $y\mapsto y+\phi_x(y)$. Its image is a complement $W_x$ of $V_x$ in $F^n$, and any complement can be obtained in this way. The nilpotent endomorphism $u'_x$ of $V'_x$ transfers to a nilpotent endomorphism $w_x$ of $W_x$.

All in all, what the given pair $(u,x)$ furnishes is a subspace $V_x$ with a basis $(x_1=x,x_2=u(x),\dots)$, a complement $W_x$, and a nilpotent endomorphism $w_x$ of $W_x$. This is more or less what the Fitting decomposition of an endomorphism gives us! Recall that $V_x$ was assumed to have been given a basis $(e_1,\dots,e_p)$. There exists a unique automorphism of $V_x$ which maps $e_i$ to $u^{i-1}(x)$ for all $i$. In other words, we have a pair of complementary subspaces $(V_x,W_x)$, a linear automorphism of $V_x$, and a nilpotent automorphism of $W_x$. By the Fitting decomposition, these data furnish in a bijective way an endomorphism of $F^n$, and that concludes the proof.

A remark about motivic integration

The framework of motivic integration suggests to upgrade these combinatorial results into equalities valid for all field $F$, which hold in the Grothendieck ring of varieties $\KVar_F$. As an abelian group, it is generated by symbols $[X]$, for all algebraic varieties $X$ over $F$, with relations $[X]=[Y]+[X\setminus Y]$, whenever $Y$ is a closed subvariety of $X$. The ring structure is defined so that the formula $[X]\cdot[Y]=[X\times Y]$ for all algebraic varieties $X$ and $Y$ over $F$.

By construction of this ring, equalities $[X]=[Y]$ in $\KVar_F$ imply that many invariants of $X$ and $Y$ coincide. In particular, when $F$ is a finite field, they will have the same number of points.

The question is thus to compute the class $[\mathscr N_n]$ of the variety $\mathscr N_n$, for any field $F$. The proofs that I described above can be more or less transferred to this context and imply the following theorem. We denote by $\mathbf L\in \KVar_F$ the class of the affine line $\mathbf A^1$.

Theorem. — One has an equality $[\mathscr N_n] \mathbf L^n = \mathbf L^{n^2}$ in the localization of the Grothendieck ring $\KVar_F$ by the element $(\mathbf L-1)\dots(\mathbf L^{n-1}-1)$.

The following question is then natural. (I have not thought about it at all.)

Question. — Does one have $[\mathscr N_n]=\mathbf L^{n^2-n}$ in $\KVar_F$?

A combinatorial proof of the Newton identities

Let $T_1,\dots,T_n$ be indeterminates. For all $m\geq0$, denote by $S_m$ the $m$th elementary symmetric polynomial, given by \[ S_m = \sum_{i_1 \lt \dots \lt i_m} T_{i_1}\dots T_{i_m}. \] These are polynomials in $T_1,\dots, T_n$, and as their names suggest, these polynomials are symmetric, meaning that \[ S_m (T_{\sigma(1)},\dots,T_{\sigma(n)}) = S_m (T_1,\dots,T_n) \] for any permutation $\sigma$ of $\{1,\dots,n\}$. By definition, one has $S_0=1$ (there is exactly one family of integers, the empty one, and the empty product is equal to~$1$) and $S_m=0$ for $m>n$ (there are no families of integers $1\leq i_1\lt\dots\lt i_m\leq n$). A well-known theorem asserts that any symmetric polynomial in $T_1,\dots,T_n$ (with coefficients in a commutative ring $R$) can be written uniquely as a polynomial in $S_1,\dots,S_n$ (with coefficients in $R$). It is in particular the case for the Newton polynomials, defined for $p\geq 0$ by \[ N_p = T_1^p + \dots+T_n^p . \] In particular, $N_0=n$ and $N_1=S_1$. The next Newton polynomial can be computed by “completing the square”: \[ N_2 = T_1^2+\dots+T_n^2 = (T_1+\dots+T_n)^2 - 2 S _1 = S_1^2 - 2S_1. \] These are the first two (or three) of a family of relations, called the Newton identities, that relate the polynomials $N_p$ and the polynomials $S_m$: \[  \sum_{m=0}^{p-1} (-1)^m N_{p-m}\cdot S_m + (-1)^p p \cdot S_p = 0. \] Using that the coefficient of $N_p$ is given by $S_0=1$, they allow to express inductively the polynomials $N_p$ in terms of $S_1,\dots,S_p$. If $p>n$, all terms with $m>n$ vanish. Using that the coefficient of $S_p$ is $N_0=p$, these formulas also show, conversely, that if $p!$ is invertible in the commutative ring $R$, then $S_1,\dots,S_p$ can be expressed in terms of $N_1,\dots,N_p$. There are various proofs of these identities. Most of the one I have seen are algebraic in nature, in so that they rely on the relations between the roots of a polynomial and its coefficients, and/or on expansion in power series. The following proof, due to Doron Zeilberger (1984) (Discrete Math. 49, p. 319, doi:10.1016/0012-365X(84)90171-7), is purely combinatorial. I have been made aware of it because it is the one that is formalized in the proof assistant system Lean. (See there for the source code.) We consider the set $S=\{1,\dots,n\}$ and define a set $X$ whose elements are the pairs $(A,k)$ satisfying the following properties:

• $A$ is a subset of $S$;
• $k$ is an element of $S$;
• the cardinality of $A$ is less or equal than $p$;
• if ${\mathop {\rm Card}}(A)=p$, then $k\in A$.

Define a map $f\colon X\to X$ by the following rule: for $(A,k)$ in $X$, set

• $f(A,k) = (A \setminus\{k\}, k)$ if $k\in A$;
• $f(A,k) = (A\cup\{k\}, k)$ if $k\notin A$.

Write $f(A,k)=(A',k)$; observe that this is again an element of $X$. The first two conditions are obvious. In the first case, when $k\in A$, then the cardinality of $A'$ is one less than the cardinality of $A$, hence is at most $p-1$, and the last condition is vacuous. In the second case, when $k\notin A$, the cardinality of $A$ is strictly smaller than $p$, because $(A,k)\in X$, so that the cardinality of $A'$ is at most $p$; moreover, $k\in A'$ hence the last condition holds. Observe that $f\colon X\to X$ is an involution: $f\circ f={\mathrm{id}}_X$. Indeed, with the same notation, one has $f(A,k)=(A',k)$. In the first case, $k\in A$, hence $k\notin A'$, so that $f(A',k)=(A'\cup\{k\},k)=(A,k)$; in the second case, one has $m\in A'$, hence $f(A',k)=(A'\setminus\{k\}, k)=(A,k)$ since $A'=A\cup\{k\}$ and $k\notin A$. Moreover, $f$ has no fixed point because it switches pairs $(A,k)$ such that $k\in A$ and pairs such that $k\notin A$. Zeilberger's proof of the Newton identities build on a function $w\colon X\to R[T_1,\dots,T_n]$ such that $w(f(A,k))=-f(A,k)$ for all $(A,k)\in X$. Since $f$ has no fixed point, the expression vanishes \[ \sum_{(A,k)\in X} w(A,k) \] since one can cancel a term $(A,k)$ with its image $f(A,k)$. Zeilberger's definition of the function $w$ is \[ w(A,k) = (-1)^{\mathrm{Card}(A)} T_k^{p-\mathrm{Card}(A)}\cdot \prod_{a\in A} T_a. \] It satisfies the desired relation: for $(A,k)\in X$ such that $k\notin A$, one has $f(A,k)=(A\cup\{k\}, k)$, so that \[ w(f(A,k)) = (-1)^{\mathrm{Card}(A)+1} T_k^{k-\mathrm{Card}(A)-1} \prod_{a\in A\cup\{k\}} T_a = - (-1)^{\mathrm{Card}(A)} T_k^{k-\mathrm{Card}(A)} \prod_{a\in A} T_a = w(A,k). \] When $k\in A$, one has $f(A,k)=(A',k)$ with $k\notin A'$, and the formula applied to $(A',k)$ implies the one for $(A,k)$, because $f$ is an involution. Now, in the relation $ \sum_{(A,k)} w(A,k) =0$, we first sum on $A$ and then on $k$: \[ \sum_{A\subset S} \sum_{k | (A,k) \in X} (-1)^{\mathrm{Card}(A)} T_k^{p-{\mathrm{Card}(A)}} \prod_{a\in A} T_a = 0\] and put together all those sets $A$ which have the same cardinality, say, $m$. This gives \[ \sum_{m=0}^n (-1)^m \sum_{\mathrm{Card}(A)=m} \prod_{a\in A} T_a \cdot \sum_{k | (A,k)\in X} T_k^{p-m} = 0. \] Let us compute separately the contribution to the sum of each $m$. When $m\gt p$, no subset $A$ is authorized so that the corresponding sum is zero. On the opposite, when $m\lt p$, for any subset $A$ such that $\mathrm{Card}(A)=m$, all pairs $(A,k)$ belong to $X$, so that the inner sum is $N_{p-m}$, and the whole term is $(-1)^mS_m N_{p-m}$. Finally, when $m=p$, the only $k$ that appear in the inner sum are the elements of $A$, and the corresponding term is $\sum_{k\in A} T_k^0=\mathrm {Card}(A)=p$; this term contributes to $(-1)^p p S_p$.

On the Hadamard-Lévy theorem, or is it Banach-Mazur?

During the preparation of an agrégation lecture on connectedness, I came across the following theorem, attributed to Hadamard–Lévy: 

Theorem. — Let $f\colon \mathbf R^n\to\mathbf R^n$ be a $\mathscr C^1$-map which is proper and a local diffeomorphism. Then $f$ is a global diffeomorphism.

In this context, that $f$ is proper means that $\| f(x)\| \to+\infty$ when $\| x\|\to+\infty$, while, by the inverse function theorem, the condition that $f$ is a local diffeomorphism is equivalent to the property that its differential $f'(x)$ is invertible, for every $x\in\mathbf R^n$. The conclusion is that $f$ is a diffeomorphism from $\mathbf R^n$ to itself; in particular, $f$ is bijective and its inverse is continuous.

This theorem is not stated in this form neither by Hadamard (1906), nor by Lévy (1920), but is essentially due to Banach & Mazur (1934) and it is the purpose of this note to clarify the history, explain a few proofs, as well as more recent consequences for partial differential equations.

A proper map is closed: the image $f(A)$ of a closed subset $A$ of $\mathbf R^n$ is closed in $\mathbf R^n$. Indeed, let $(a_m)$ be a sequence in $A$ whose image $(f(a_m))$ converges in $\mathbf R^n$ to an element $b$; let us show that there exists $a\in A$ such that $b=f(a)$. The properness assumption on $f$ implies that $(a_m)$ is bounded. Consequently, it has a limit point $a$, and $a\in A$ because $A$ is closed. Necessarily, $f(a)$ is a limit point of the sequence $(f(a_m))$, hence $b=f(a)$.

In this respect, let us note the following reinforcement of the previous theorem, due to Browder (1954):
Theorem (Browder). — Let $f\colon \mathbf R^n\to\mathbf R^n$ be a local homeomorphism. If $f$ is closed, then $f$ is a global homeomorphism.

A surprising aspect of these results and their descendents is that they are based on two really different ideas. Banach & Mazur and Browder are based on the notion of covering, with ideas of homotopy theory and, ultimately, the fact that $\mathbf R^n$ is simply connected. On the other hand, the motivation of Hadamard was to generalize to dimension $n$ the following elementary discussion in the one-dimensional case: Let $f\colon\mathbf R\to\mathbf R$ be a $\mathscr C^1$-function whose derivative is $>0$ everywhere (so that $f$ is strictly increasing); give a condition for $f$ to be surjective. In this case, the condition is easy to find: the indefinite integral $\int f'(x)\,dx$ has to be divergent both at $-\infty$ and $+\infty$. In the $n$-dimensional case, the theorems of Hadamard is the following:

Theorem.Let $f\colon\mathbf R^n\to\mathbf R^n$ be a $\mathscr C^1$-map. For $r\in\mathbf R_+$, let $\omega(r)$ be the infimum, for $x\in\mathbf R^n$ such that $\|x\|=r$, of the norm of the linear map $f'(x)^{-1}$; if $\int_0^\infty dr/\omega(r)=+\infty$, then $f$ is a global diffeomorphism.

In Hadamard's paper, the quantity $\omega(r)$ is described geometrically as the minor axis of the ellipsoid defined by $f'(x)$, and Hadamard insists that using the volume of this ellipsoid only, essentially given by the determinant of $f'(x)$, would not suffice to characterize global diffeomorphisms. (Examples are furnished by maps of the form $f(x_1,x_2)=(f_1(x_1),f_2(x_2))$. The determinant condition considers $f_1'(x_1)f_2'(x_2)$, while one needs individual conditions on $f'_1(x_1)$ and $f'_2(x_2)$.)

In fact, as explained in Plastock (1974), both versions (closedness hypothesis or quantitative assumptions on the differential) imply that the map $f$ is a topological covering of $\mathbf R^n$. Since the target $\mathbf R^n$ is simply connected and the source $\mathbf R^n$ is connceted, $f$ has to be a homeomorphism. I will explain this proof below, but I would first like to explain another one, due to Zuily & Queffelec (1995) propose an alternating proof which is quite interesting.

A dynamical system approach

The goal is to prove that $f$ is bijective and, to that aim, we will prove that every preimage set $f^{-1}(b)$ is reduced to one element. Replacing $f$ by $f-b$, it suffices to treat the case of $b=0$. In other words, we wish to solve that the equation $f(x)=0$ has exactly one solution. For that, it is natural to try to start from some point $\xi\in\mathbf R^n$ and to force $f$ to decrease. This can be done by following the flow of the vector field given by $v(x)=-f'(x)^{-1}(f(x))$. This is a vector field on $\mathbf R^n$ and we can consider its flow: a map $\Phi$ defined on an open subset of $\mathbf R\times\mathbf R^n$ such that $\partial_t \Phi(t,x)=v(\Phi(t,x))$ for all $(t,x)$ and $\Phi(0,x)=x$ for all $x$. In fact, the Cauchy–Lipschitz theorem guarantees the existence of such a flow only if the vector field $v$ is locally Lipschitz, which happens if, for example, $f$ is assumed to be $\mathscr C^2$. In this case, there is even uniqueness of a maximal flow, and we will make this assumption, for safety. (In fact, the paper of De Marco, Gorni & Zampieri (1994) constructs the flow directly thanks to the hypothesis that the vector field is pulled back from the Euler vector field on $\mathbf R^n$.)

What are we doing here? Note that in $\mathbf R^n$, the opposite of the Euler vector field, defined by $u(y)=-y$, has a very simple solution: the flow lines are straight lines going to $0$. The formula above just pulls back this vector field $u$ via the local diffeomorphism $f$, and the flow lines of the vector field $v$ will just be the ones given by pull back by $f$, which will explain the behaviour described below.

In particular, let $a\in\mathbf R^n$ be such that $f(a)=0$ and let $U$ be a neighborhood of $a$ such that $f$ induces a diffeomorphism from $U$ to a ball around $0$. Pulling back the solution of the minus-Euler vector field by $f$, we see that once a flow line enters the open set $U$, it converges to $a$. The goal is now to prove that it will indeed enter such a neighborhood (and, in particular, that such a point $a$ exists).

We consider a flow line starting from a point $x$, that is, $\phi(t)=\Phi(t,x)$ for all times $t$. Let $g(t)= f(\phi(t))$; observe that $g$ satisfies $g'(t)=f'(\phi(t))(\phi'(t))=-g(t)$, hence $g(t)=g(0)e^{-t}$. Assume that the line flow is defined on $[0;t_1\mathopen[$, with $t_1<+\infty$. by what precedes, $g$ is bounded in the neighborhood of $t_1$; since $f$ is assumed to be proper, this implies that $\phi(t)$ is bounded as well. The continuity of the vector field $v$ implies that $\phi$ is uniformly continuous, hence it has a limit at $t_1$. We may then extend the line flow a bit right of $t_1$. As a consequence, the line flow is defined for all times, and $g(t)\to0$ when $t\to+\infty$. By the same properness argument, this implies that $\phi(t)$ is bounded when $t\to+\infty$, hence it has limit points $a$ which satisfy $f(a)=0$. Once $\phi$ enters an appropriate neighborhood of such a point, we have seen that the line flow automatically converges to some point $a\in f^{-1}(0)$.

Let us now consider the map $\lambda\colon\mathbf R^n\to f^{-1}(0)$ that associates with a point $\xi$ the limit of the line flow $t\mapsto \Phi(t,\xi)$ starting from the initial condition $\xi$. By continuity of the flow of a vector field depending on the initial condition, the map $\lambda$ is continuous. On the other hand, the hypothesis that $f$ is a local diffeomorphism implies that $f^{-1}(0)$ is a closed discrete subset of $\mathbf R^n$. Since $\mathbf R^n$ is connected, the map $\lambda$ is constant. Since one has $\lambda(\xi)=\xi$ for every $\xi\in f^{-1}(0)$, this establishes that $f^{-1}(0)$ is reduced to one element, as claimed.

Once $f$ is shown to be bijective, the fact that it is proper (closed would suffice) implies that its inverse bijection $f^{-1}$ is continuous. This concludes the proof.

The theorem of Banach and Mazur

The paper of Banach and Mazur is written in a bigger generality. They consider multivalued continuous maps $F\colon X\to Y$ ($k$-deutige stetige Abbildungen) by which they mean that for every $x$, a subset $F(x)$ of $Y$ is given, of cardinality $k$, the continuity being expressed by sequences: if $x_n\to x$, one can order, for every $n$, the elements of $F(x_n)=\{y_{n,1},\dots,y_{n,k}\}$, as well as the elements of $F(x)=\{y_1,\dots,y_k\}$, in such a way that $y_{n,j}\to y_n$ for all $j$. (In their framework, $X$ and $Y$ are metric spaces, but one could transpose their definition to topological spaces if needed.) They say that such a map is decomposed (zerfällt) if there are continuous functions $f_1,\dots,f_k$ from $X$ to $Y$ such that $F(x)=\{f_1(x),\dots,f_k(x)\}$ for all $x\in X$.

In essence, the definition that Banach and Mazur are proposing contains as a particular case the finite coverings. Namely, if $p\colon Y\to X$ is a finite covering of degree $k$, then the map $x\mapsto p^{-1}(x)$ is a continuous $k$-valued map from $X$ to $Y$. Conversely, let us consider the graph $Z$ of $F$, namely the set of all points $(x,y)\in X\times Y$ such that $y\in F(x)$. Then the first projection $p\colon Z\to X$ is a covering map of degree $k$, but it is not clear that it has local sections.

It would however not be so surprising to 21st-century mathematicians that if one makes a suitable assumption of simple connectedness on $X$, then every such $F$ should be decomposed. Banach and Mazur assume that $X$ satisfies two properties:

1. The space $X$ is semilocally arcwise connected: for every point $x\in X$ and every neighborhood $U$ of $x$, there exists an open neighborhood $U'$ contained in $U$ such that for every point $x'\in U'$, there exists a path $c\colon[0;1]\to U$ such that $c(0)=x$ and $c(1)=x'$. (Semilocally means that the path is not necessarily in $U'$ but in $U$.)
2. The space $X$ is arcwise simply connected: two paths $c_0,c_1\colon[0;1]\to X$ with the same endpoints ($c_0(0)=c_1(0)$ and $c_0(1)=c_1(1)$) are strictly homotopic — there exists a continuous map $h\colon[0;1]\to X$ such that $h(0,t)=c_0(t)$ and $h(1,t)=c_1(t)$ for all $t$, and $h(s,0)=c_0(0)$ and $h(s,1)=c_0(1)$ for all $s$.

Consider a $k$-valued continuous map $F$ from $X$ to $Y$, where $X$ is connected. Banach and Mazur first prove that for every path $c\colon [0;1]\to X$ and every point $y_0\in F(c(0))$, there exists a continuous function $f\colon[0;1]\to Y$ such that $f(t)\in F(c(t))$ for all $t$. To that aim, the consider disjoint neighborhoods $V_1,\dots,V_k$ of the elements of $F(c(0))$, with $y_0\in V_1$, say, and observe that for $t$ small enough, there is a unique element in $F(c(t))\cap V_1$. This defines a bit of the path $c$, and one can go on. Now, given two paths $c,c'$ such that $c(0)=c'(0)$ and $c(1)=c'(1)$, and two maps $f,f'$ as above, they consider a homotopy $h\colon[0;1]\times[0;1]\to X$ linking $c$ to $c'$. Subdividing this square in small enough subsquares, one see by induction that $f(1)=f'(1)$. (This is analogous to the proof that a topological covering of the square is trivial.) Fixing a point $x_0\in X$ and a point $y_0\in F(x_0)$, one gets in this way a map from $X$ to $Y$ such that $F(x)$ is equal to $f(1)$, for every path $c\colon[0;1]\to X$ such that $c(0)=x_0$ and $c(1)=x$, and every continuous map $f\colon [0;1]\to Y$ such that $f(t)\in F(c(t))$ for all $t$ and $f(0)=y_0$. This furnishes a map from $X$ to $Y$, and one proves that it is continuous. If one considers all such maps, for all points in $F(x_0)$, one obtains the decomposition of the multivalued map $F$.

To prove their version of the Hadamard–Lévy theorem, Banach and Mazur observe that if $f\colon Y\to X$ is a local homeomorphism which is proper, then setting $F(x)=f^{-1}(y)$ gives a multivalued continuous map. It is not obvious that the cardinalities $k(x)$ of the sets $F(x)$ are constant, but this follows (if $X$ is connected) from the fact that $f$ is both a local homeomorphism and proper. Then $F$ is decomposed, so that there exist continuous maps $g_1,\dots,g_k\colon X\to Y$ such that $f^{-1}(x)=\{g_1(x),\dots,g_k(x)\}$ for all $x\in X$. This implies that $Y$ is the disjoint union of the $k$ connected subsets $g_j(X)$. If $Y$ is connected, then $f$ is a homeomorphism.

The versions of Hadamard and Lévy, after Plastock

Hadamard considered the finite dimensional case, and Lévy extended it to the case of Hilbert spaces.

Plastock considers a Banach-space version of the theorem above: $f\colon E\to F$ is a $\mathscr C^1$-map between Banach spaces with invertible differentials and such that, setting $\omega(r)=\inf_{\|x\| = r}\|f'(x)^{-1}\|$, one has $\int_0^\infty \omega(r)\,dr=+\infty$. Of course, under these hypotheses, the Banach spaces $E$ and $F$ are isomorphic, but it may be useful that they are not identical. Note that $f(E)$ is open in $F$, and the proposition that will insure that $f$ is a global diffeomorphism is the following one, in the spirit of covering theory.

Proposition. – (Assuming that $f$ is a local diffeomorphism.) It suffices to prove that the map $f$ satisfies the path lifting property: for every point $x\in E$ and every $\mathscr C^1$ map $c\colon[0;1]\to f(E)$ such that $c(0)=f(x)$, there exists a $\mathscr C^1$ map $d\colon[0;1]\to E$ such that $c(t)=f(d(t))$ for all $t$ and $d(0)=c$.

The goal is now to prove that $f$ satisfies this path lifting property. Using that $f$ is a local homeomorphism, one sees that lifts are unique, and are defined on a maximal subinterval of $[0;1]$ which is either $[0;1]$ itself, or of the form $[0;s\mathclose[$. To prevent the latter case, one needs to impose conditions on the norm $\| f'(x)^{-1}\|$ such as the one phrased in terms of $\omega(r)$ as in the Hadamard–Lévy theorem. In fact, Plastock starts with a simpler case.

Proposition. – The path lifting property follows from the following additional hypotheses:

1. One has $\|f(x)\|\to+\infty$ when $\|x\|\to+\infty$;
2. There exists a positive continuous function $M\colon\mathbf R_+\to\mathbf R_+$ such that $\|f'(x)^{-1}\|\leq M(\|x\|)$ for all $x.

Assume indeed that a path $c$ has a maximal lift $d$, defined over the interval $[0;s\mathclose[$. By the hypothesis (i), $d(t)$ remains bounded when $t\to s$, because $c(t)=f(d(t))$ tends to $c(s)$. Differentiating the relation $c(t)=f(d(t))$, one gets $c'(t)=f'(d(t))(d'(t))$, hence $d'(t)=f'(d(t))^{-1}(c'(t))$, so that $\| d'(t)\|\leq M(\|d(t)\|) \|c'(t)\|$. This implies that $\|d'\|$ is bounded, so that $d$ is uniformly continuous, hence it has a limit at $s$. Then the path $d$ can be extended by setting $d(s)$ to this limit and using the local diffeomorphism property to go beyong $s$.

The Hadamard–Lévy is related to completeness of some length-spaces. So we shall modify the distance of the Banach space $E$ as follows: if $c\colon[0;1]\to E$ is a path in $E$, then its length is defined by \[ \ell(c) = \int_0^1 \| f'(c(t))^{-1}\|^{-1} \|{c'(t)}\|\, dt. \] Observe that $\|f'(c(t))^{-1}\|^{-1} \geq \omega(\|c(t)\|)$, so that \[ \ell(c) \geq \int_0^1 \omega(\|c(t)\|) \|{c'(t)}\|\, dt. \] The modified distance of two points in $E$ is then redefined as the infimum of the lengths of all paths joining two points.

Lemma. – With respect to the modified distance, the space $E$ is complete.

One proves that $\ell(c) \geq \int_{\|{c(0)}\|}^{\|{c(1)}\|}\omega(r)\,dr$. Since $\int_0^\infty \omega(r)\,dr=+\infty$, this implies that Cauchy sequences for the modified distance are bounded in $E$ for the original norm. On the other hand, on any bounded subset of $E$, the Banach norm and the modified distance are equivalent, so that they have the same Cauchy sequences.

Other conditions can be derived from Plastock's general theorem. For example, assuming that $E$ and $F$ are a Hilbert space $H$, he shows that it suffices to assume the existence of a decreasing function $\lambda\colon\mathbf R_+\to\mathbf R_+$ such that $\langle f'(x)(u),u\rangle \geq \lambda(\|x\|) \| u\|^2$ for all $x,y$ and $\int_0^\infty \lambda(r)\,dr=+\infty$. Indeed, under this assumption, one may set $\omega(r)=\lambda(r)$.

Application to periodic solutions of differential equations

Spectral theory can be seen as the infinite dimensional generalization of classical linear algebra. Linear differential operators and linear partial differential operators furnish prominent examples of such operators. The theorems of Hadamard–Lévy type have been applied to solve nonlinear differential equations.

I just give an example here, to give an idea of how this works, and also because I am quite lazy enough to check the details.

Following Brown & Lin (1979), we consider the Newtonian equation of motion: \[ u''(t) + \nabla G (u(t)) = p(t) \] where $G$ represents the ambiant potential, assumed to be smooth enough, and $p\colon \mathbf R\to\mathbf R^n$ is some external control. The problem studied by Brown and Lin is to prove the existence of periodic solutions when $p$ is itself periodic. The method consists in interpreting the left hand side as a non linear map defined on the Sobolev space $E$ of $2\pi$-periodic $\mathscr C^1$-functions with a second derivative in $F=L^2([0;2\pi];\mathbf R^n)$, with values in $F$. Write $L$ for the linear operator $u\mapsto u''$ and $N$ for the (nonlinear) operator $u\mapsto \nabla G(u)$. Then $L$ is linear continuous (hence $L'(u)(v)=L'(v)$), and $N$ is continuously differentiable, with differential given by \[ N'(u) (v) = \left( t \mapsto Q (u(t)) (v(t)) \right) \] for $u,v\in E$, and $Q$ is the Hessian of $G$.

In other words, the differential $(L+N)'(u)$ is the linear map $v\mapsto L(v) + Q(u(t)) v$. It is invertible if the eigenvalues of $Q(u(t))$ are away from integers. Concretely, Brown and Lin assume that there are two constant symmetric matrices $A$ and $B$ such that $A\leq Q(x) \leq B$ for all $x$, and whose eigenvalues $\lambda_1\leq \dots\lambda_n$ and $\mu_1\leq\dots\leq \mu_n$ are such that there are integers $N_1,\dots,N_n$ with $N_k^2<\lambda_k\leq\mu_k<(N_k+1)^2$ for all $k$. Using spectral theory in Hilbert spaces, these conditions imply that the linear operator $L+Q(u)\colon E\to F$ is an isomorphism, and that $\|(L+Q(u)^{-1}\|$ is bounded from above by the constant expression \[ c= \sup_{1\leq k\leq n} \sup (\lambda_k-N_k^2)^{-1},((N_k+1)^2-\mu_k)^{-1} ).\]

Thanks to this differential estimate, the theorem of Hadamard–Lévy implies that the nonlinear differential operator $L+N$ is a global diffeomorphism from $E$ to $F$. In particular, there is a unique $2\pi$-periodic solution for every $2\pi$-periodic control function $p$.

I thank Thomas Richard for his comments.

On Rolle's theorem

This post is inspired by a paper of Azé and Hiriart-Urruty published in a French high school math journal; in fact, it is mostly a paraphrase of that paper with the hope that it be of some interest to young university students, or to students preparing Agrégation. The topic is Rolle's theorem.

1. The one-dimensional theorem, a generalization and two other proofs

Let us first quote the theorem, in a nonstandard form.

Theorem. — Let $I=\mathopen]a;b\mathclose[$ be a nonempty but possibly unbounded interval of $\mathbf R$ and let $f\colon I\to\mathbf R$ be a continuous function. Assume that $f$ has limits at $a$ and $b$, equal to some element $\ell\in\mathbf R\cup\{+\infty\}$. Then $f$ is bounded from below.

1. If $\inf_I(f)<\ell$, then there exists a point $c\in I$ such that $f(c)=\inf_I (f)$. If, moreover, $f$ has a right derivative and a left derivative at $c$, then $f'_l(c)\leq0$ and $f'_r(c)\geq0$.
2. If $\inf_I(f)\geq\ell$, then $f$ is bounded on $I$ and there exists a point $c\in I$ such that $f(c)=\sup_I(f)$. If, moreover, $f$ has a right derivative and a left derivative at $c$, then $f'_l(c)\geq0$ and $f'_r(c)\leq0$.

Three ingredients make this version slightly nonstandard:

• The interval $I$ may be taken to be infinite;
• The function $f$ may tend to $+\infty$ at the endpoints of $I$;
• Only left and  right derivatives are assumed.

Of course, if $f$ has a derivative at each point, then the statement implies that $f'(c)=f'_l(c)=f'_r(c)=0$.

a) As stated in this way, the proof is however quite standard and proceeds in two steps.

1. Using that $f$ has a limit $\ell$ which is not $-\infty$ at $a$ and $b$, it follows that there exists $a'$ and $b'$ in $I$ such that $a<a'<b'<b$ such that $f$ is bounded from below on $\mathopen ]a;a']$ and on $[b';b\mathclose[$. Since $f$ is continuous on the compact interval $[a';b']$, it is then bounded from below on $I$.
   If $\inf_I(f)<\ell$, then we can choose $\ell'\in\mathbf R$ such that $\inf_I(f)<\ell'<\ell$ and $a'$, $b'$ such that $f(x)>\ell'$ outside of $[a';b']$. Then, let $c\in [a';b']$ such that $f(c)=\inf_{[a';b']}(f)$; then $f(c)=\inf_I(f)$. 
   If $\sup_I(f)>\ell$, then we have in particular $\ell\neq+\infty$, and we apply the preceding analysis to $-f$.
   In the remaining case, $\inf_I(f)=\sup_I(f)=\ell$ and $f$ is constant.
   
2. For $x>c$, one has $f(x)\geq f(c)$, hence $f'_r(c)\geq 0$; for $x<c$, one has $f(x)\geq f(c)$, hence $f'_l(c)\leq0$.

The interest of the given formulation can be understood by looking at the following two examples.

1. If $f(x)=|x|$, on $\mathbf R$, then $f$ attains its lower bound at $x=0$ only, where one has $f'_r(0)=1$ and $f'_l(0)=-1$.
2. Take $f(x)=e^{-x^2}$. Then there exists $c\in\mathbf R$ such that $f'(c)=0$. Of course, one has $f'(x)=-2xe^{-x^2}$, so that $c=0$. However, it is readily seen by induction that for any integer $n$, the $n$th derivative of $f$ is of the form $P_n(x)e^{-x^2}$, where $P_n$ has degree $n$. In particular, $f^{(n)}$ tends to $0$ at infinity. And, by induction again, the theorem implies that $P_n$ has $n$ distinct roots in $\mathbf R$, one between any two consecutive roots of $P_{n-1}$, one larger than the largest root of $P_n$, and one smaller than the smallest root of $P_n$.

b) In a 1959 paper, the Rumanian mathematician Pompeiu proposed an alternative proof of Rolle's theorem, when the interval $I$ is bounded, and which works completely differently. Here is how it works, following the 1979 paper published in American Math. Monthly by Hans Samelson.

First of all, one uses the particular case $n=2$ of the Levi chord lemma :

Lemma. — Let $f\colon [a;b]\to\mathbf R$ be a continuous function such that $f(a)=f(b)$. For every integer $n\geq 2$, there exists $a',b'\in[a;b]$ such that $f(a')=f(b')$ and $b'-a'=(b-a)/n$.

Let $h=(b-a)/n$. From the equality
\[ 0 = f(b)-f(a) = (f(a+h)-f(a))+(f(a+2h)-f(a+h))+\cdots + (f(a+nh)-f(a+(n-1)h), \]
one sees that the function $x\mapsto f(x+h)-f(x)$ from $[a;b-h]$ to $\mathbf R$ does not have constant sign. By the intermediate value theorem, it vanishes at some point $a'\in [a;b-h]$. If $b'=a'+h$, then $b'\in[a;b]$, $b'-a'=(b-a)/n$ and $f(a')=f(b')$.

Then, it follows by induction that there exists a sequence of nested intervals $([a_n;b_n])$ in $[a;b]$ with $f(a_n)=f(b_n)$ and $b_n-a_n=(b-a)/2^n$ for all $n$. The sequences $(a_n)$ and $(b_n)$ converge to a same limit $c\in [a;b]$. Since $f(b_n)=f(c)+(b_n-c) (f'(c) + \mathrm o(1))$, $f(a_n)=f(c)+(a_n-c)(f'(c)+\mathrm o(1))$, one has
\[ f'(c) = \lim \frac{f(b_n)-f(a_n)}{b_n-a_n} = 0. \]

What makes this proof genuinely distinct from the classical one is that the obtained point $c$ may not be a local minimum or maximum of $f$, also I don't have an example to offer now.

c) In 1979, Abian furnished yet another proof, which he termed as the “ultimate” one. Here it is:

It focuses on functions $f\colon[a;b]\to\mathbf R$ on a bounded interval of $\mathbf R$ which are not monotone and, precisely, which are up-down, in the sense that $f(a)\leq f(c)$ and $f(c)\geq f(b)$, where $c=(a+b)/2$ is the midpoint of $f$. If $f(a)=f(b)$, then either $f$ or $-f$ is up-down.

Then divide the interval $[a;b]$ in four equal parts: $[a;p]$, $[p;c]$, $[c;q]$ and $[q;b]$. If $f(p)\geq f(c)$, the $f|_{[a;c]}$ is up-down. Otherwise, one has $f(p)\leq f(c)$. In this case, if $f(c)\geq f(q)$, we see that $f|_{[p;q]}$ is up-down. And otherwise, we observe that $f(q)\leq f(c)$ and $f(c)\geq f(b)$, so that $f|_{[c;b]}$ is up-down. Conclusion: we have isolated within the interval $[a;b]$ a subinterval $[a';b']$ of length $(b-a)/2$ such that $f|_{[a';b']}$ is still up-down.

Iterating the procedure, we construct a sequence $([a_n;b_n])$ of nested intervals, with $(b_n-a_n)=(b-a)/2^n$ such that the restriction of $f$ to each of them is up-down. Set $c_n=(a_n+b_n)/2$.

The sequences $(a_n), (b_n),(c_n)$ satisfy have a common limit $c\in [a;b]$. From the inequalities $f(a_n)\leq f(c_n)$ and $a_n\leq c_n$,  we obtain $f'(c)\geq 0$; from the inequalities $f(c_n)\geq f(b_n)$ and $c_n\leq b_n$, we obtain $f'(c)\leq 0$. In conclusion, $f'(c)=0$.

2. Rolle's theorem in normed vector spaces

Theorem. — Let $E$ be a normed vector space, let $U$ be an open subset of $E$ and let $f\colon U\to\mathbf R$ be a differentiable function. Assume that there exists $\ell\in\mathbf R\cup\{+\infty\}$ such that $f(x)\to \ell$ when $x$ tends to the “boundary” of $U$ — for every $\ell'<\ell$, there exists a compact subset $K$ of $U$ such that $f(x)\geq\ell'$ for all $x\in U$ but $x\not\in K$. Then $f$ is bounded below on $U$, there exists $a\in U$ such that $f(a)=\inf_U (f)$ and $Df(a)=0$.

The proof is essentially the same as the one we gave in dimension 1. I skip it here.

If $E$ is finite dimensional, then this theorem applies in a vast class of examples : for example, bounded open subsets $U$ of $E$, and continuous functions $f\colon \overline U\to\mathbf R$ which are constant on the boundary $\partial(U)=\overline U - U$ of $U$ and differentiable on $U$.

However, if $E$ is infinite dimensional, the closure of a bounded open set is no more compact, and it does not suffice that $f$ extends to a function on $\overline U$ with a constant value on the boundary.

Example. — Let $E$ be an infinite dimensional Hilbert space, let $U$ be the open unit ball and $B$ be the closed unit ball. Let $g(x)=\frac12 \langle Ax,x\rangle+\langle b,x\rangle +c$ be a quadratic function, where $A\in\mathcal L(E)$, $b\in E$ and $c\in\mathbf R$, and let $f(x)=(1-\lVert x\rVert^2) g(x)$. The function $f$ is differentiable on $E$ and one has
\[  \nabla f(x) =  (1-\Vert x\rVert^2) ( Ax + b) - 2 (\frac12 \langle Ax,x\rangle + \langle b,x\rangle + c) x. \]
Assume that there exists $x\in U$ such that $\nabla f(x)=0$. Then $Ax+b = \lambda x$, with
\[ \lambda= \frac2{1-\lVert x\rVert ^2} \left(\frac12 \langle Ax,x\rangle + \langle b,x\rangle + c \right). \]
Azé and Hiriart-Urruty take $E=L^2([0;1])$, for $A$ the operator of multiplication by the function $t$,  $b(t)=t(1-t)$, and $c=4/27$. Then, one has $g(x)>0$, hence $\lambda>0$, and $x(t)=\frac1{\lambda-t}b(t)$ for $t\in[0;1]$. This implies that $\lambda\geq 1$, for, otherwise, the function $x(t)$ would not belong to $E$. This allows to compute $\lambda$ in terms of $\mu$,  obtaining $\lambda\leq3/4$, which contradicts the inequality $\lambda\geq 1$. (I refer to the paper of Azé and Hiriart-Urruty for more details.)

3. An approximate version of Rolle's theorem

Theorem. — Let $B$ the closed euclidean unit ball in $\mathbf R^n$, let $U$ be its interior, let $f\colon B\to \mathbf R$ be a continuous function on $B$. Assume that $\lvert f\rvert \leq \epsilon $ on the boundary $\partial(U)$ and that $f$ is differentiable on $U$. Then there exists $x\in U$ such that $\lVert Df(x)\rVert\leq\epsilon$.

In fact, replacing $f$ by $f/\epsilon$, one sees that it suffices to treat the case $\epsilon =1$.

Let $g(x)=\lVert x\rVert^2- f(x)^2$. This is a continuous function on $B$; it is differentiable on $U$, with $ \nabla g(x)=2(x-f(x)\nabla f(x))$. Let $\mu=\inf_B(g)$. Since $g(0)=-f(0)^2\leq0$, one has $\mu\leq 0$. We distinguish two cases:

1. If $\mu=0$, then $\rvert f(x)\lvert \leq \lVert x\rVert$ for all $x\in B$. This implies that $\lVert\nabla f(0)\rVert\leq1$.
2. If $\mu<0$, let $x\in B$ be such that $ g(x)=\mu$; in particular, $f(x)^2\geq \lVert x\rVert^2-\mu>0$, which implies that $f(x)\neq0$. Since $g\geq0$ on $\partial(U)$, we have $x\in B$, hence $\nabla g(x)=0$. Then $x=f(x)\nabla f(x)$, hence $\nabla f(x)=x/f(x)$. Consequently,
   \[ \lVert \nabla f(x)\rVert \leq \frac{\lVert x\rVert}{f(x)}\leq \frac{\lVert x\rVert}{(\lVert x\rVert^2-\mu)^{1/2}}<1.\]

This concludes the proof. 

Thanks to the Twitter users @AntoineTeutsch, @paulbroussous and @apauthie for having indicated me some misprints and incorrections.

Warning! — Theorems ahead

Don't worry, no danger ahead! — This is just a short post about the German mathematician Ewald Warning and the theorems that bear his name.

It seems that Ewald Warning's name will be forever linked with that of Chevalley, for the Chevalley-Warning theorem is one of the rare modern results that can be taught to undergraduate students; in France, it is especially famous at the Agrégation level. (Warning published a second paper, in 1959, about the axioms of plane geometry.)

Warning's paper, Bemerkung zur vorstehenden Arbeit von Herrn Chevalley (About a previous work of Mr Chevalley), has been published in 1935 in the Publications of the mathematical seminar of Hamburg University (Abhandlungen aus dem Mathematischen Seminar der Universität Hamburg), just after the mentioned paper of Chevalley. Emil Artin had a position in Hamburg at that time, which probably made the seminar very attractive; as a matter of fact, the same 1935 volume features a paper of Weil about Riemann-Roch, one of Burau about braids, one of Élie Cartan about homogeneous spaces, one of Santalo on geometric measure theory, etc.

1. The classic statement of the Chevalley—Warning theorem is the following.

Theorem 1. – Let $p$ be a prime number, let $q$ be a power of $p$ and let $F$ be a field with $q$ elements. Let $f_1,\dots,f_m$ be polynomials in $n$ variables and coefficients in $F$, of degrees $d_1,\dots,d_m$; let $d=d_1+\dots+d_m$. Let $Z=Z(f_1,\dots,f_m)$ be their zero-set in $F^n$. If $d<n$, then $p$ divides $\mathop{\rm Card}(Z)$.

This is really a theorem of Warning, and Chevalley's theorem was the weaker consequence that if $Z\neq\emptyset$, then $Z$ contains at least two points. (In fact, Chevalley only considers the case $q=p$, but his proof extends readily.) The motivation of Chevalley lied in the possibility to apply this remark to the reduced norm of a possibly noncommutative finite field (a polynomial of degree $d$ in $d^2$ variables which vanishes exactly at the origin), thus providing a proof of Wedderburn's theorem.

a) Chevalley's proof begins with a remark. For any polynomial $f\in F[T_1,\ldots,T_n]$, let $f^*$ be the polynomial obtained by replacing iteratively $X_i^q$ by $X_i$ in $f$, until the degree of $f$ in each variable is $<q$. For all $a\in F^n$, one has $f(a)=f^*(a)$; moreover, using the fact that a polynomial in one variable of degree $<q$ has at most $q$ roots, one proves that if $f(a)=0$ for all $a\in F^n$, then $f^*=0$.

Assume now that $Z$ contains exactly one point, say $a\in F^n$, let $f=\prod (1-f_j^{q-1})$, let $g_a=\prod (1-(x_i-a_i)^{q-1})$. Both polynomials take the value $1$ at $x=a$, and $0$ elsewhere; moreover, $g_a$ is reduced. Consequently, $f^*=g$. Then
$$ (q-1)n=\deg(g_a)=\deg(f^*)\leq \deg(f)=(q-1)\sum \deg(f_j)=(q-1)d, $$
contradicting the hypothesis that $d<n$.

b) Warning's proof is genuinely different. He first defines, for any subset $A$ of $F^n$ a reduced polynomial $g_A=\sum_{a\in A} g_a=\sum_{a\in A}\prod (1-(x_i-a_i)^{q-1})$, and observes that $g_A(a)=1$ if $a\in A$, and $g_A(a)=0$ otherwise.
Take $A=Z$, so that $f^*=g_Z$. Using that $\deg(f^*)\leq \deg(f)=(q-1)d$ and the expansion
$ (x-a)^{q-1} = \sum_{i=0}^{q-1} x^i a^{q-1-i}$, Warning derives from the equality $f^*=g_Z$
the relations
\[ \sum_{a\in Z} a_1^{\nu_1}\dots a_n^{\nu_n}=0, \]
for all $(\nu_1,\dots,\nu_n)$ such that $0\leq \nu_i\leq q-1$ and $\sum\nu_i <(q-1)(n-d)$. The particular case $\nu=0$ implies that $p$ divides $\mathop{\rm Card}(Z)$. More generally:

Proposition 2. — For every polynomial $\phi\in F[T_1,\dots,T_n]$ of reduced degree $<(q-1)(n-d)$, one has $\sum_{a\in A} \phi(a) = 0$.

c) The classic proof of that result is even easier. Let us recall it swiftly. First of all, for every integer $\nu$ such that $0\leq \nu <q$, one has $\sum_{a\in F} a^\nu=0$. This can be proved in many ways, for example by using the fact that the multiplicative group of $F$ is cyclic; on the other hand, for every nonzero element $t$ of $F$, the change of variables $a=tb$ leaves this sum both unchanged and multiplied by $t^\nu,$ so that taking $t$ such that $t^\nu\neq 1$, one sees that this sum vanishes. It follows from that that for every polynomial $f\in F[T_1,\dots,T_n]$ whose degree in some variable is $\lt q-1$, one has $\sum_{a\in F^n} f(a)=0$. This holds in particular if the total degree of $f$ is $\lt (q-1)n$.
Taking $f$ as above proves theorem 1.

2. On the other hand there is a second Warning theorem, which seems to be absolutely neglected in France. It says the following:

Theorem 3. — Keep the same notation as in theorem 1. If $Z$ is nonempty, then $\mathop{\rm Card}(Z)\geq q^{n-d}$.

To prove this result, Warning starts from the following proposition:

Proposition 4. – Let $L,L'$ be two parallel subspaces of dimension $d$ in $F^n$. Then $\mathop{\rm Card}(Z\cap L)$ and $\mathop{\rm Card}(Z\cap L')$ are congruent modulo $p$.

Let $r=n-d$. Up to a change of coordinates, one may assume that $L=\{x_1=\dots=x_{r}=0\}$ and $L'=\{x_1-1=x_2=\dots=x_{r}=0\}$. Let
\[ \phi = \frac{1-x_1^{q-1}}{1-x_1} (1-x_2^{q-1})\cdots (1-x_{r}^{q-1}).\]
This is a polynomial of total degree is $(q-1)r-1<(q-1)(n-d)$. For $a\in F^n$, one has $\phi(a)=1$ if $a\in L$, $\phi(a)=-1$ if $a\in L'$, and $\phi(a)=0$ otherwise. Proposition 4 thus follows from proposition 2. It is now very easy to prove theorem 3 in the particular case where there exists one subspace $L$ of dimension $d$ such that $\mathop{\rm Card}(Z\cap L)\not\equiv 0\pmod p$. Indeed, by proposition 4, the same congruence will hold for every translate $L'$ of $L$. In particular, $\mathop{\rm Card}(Z\cap L')\neq0$ for every translate $L'$ of $L$, and there are $q^{n-d}$ distinct translates.

To prove the general case, let us choose a subspace $M$ of $F^n$ of dimension $s\leq d$ such that $\mathop{\rm Card}(Z\cap M)\not\equiv 0\pmod p$, and let us assume that $s$ is maximal.
Assume that $s< d$. Let $t\in\{1,\dots,p-1\}$ be the integer such that $\mathop{\rm Card}(Z\cap M)\equiv t\pmod p$. For every $(s+1)$-dimensional subspace $L$ of $F^n$ that contains $M$, one has $\mathop{\rm Card}(Z\cap L)\equiv 0\pmod p$, by maximaility of $s$, so that $Z\cap (L\setminus M)$ contains at least $p-t$ points. Since these subspaces $L$ are in 1-1 correspondence with the lines of the quotient space $F^n/M$, their number is equal to $(q^{n-s}-1)/(q-1)$. Consequently, \[ \mathop{\rm Card}(Z) = \mathop{\rm Card}(Z\cap M) + \sum_L \mathop{\rm Card}(Z\cap (L\setminus M)) \geq t + (p-t) \frac{q^{n-s}-1}{q-1} \geq q^{n-s-1}\geq q^{n-d}, \] as was to be shown.

3. Classic theorems seem to an everlasting source of food for thought.

a) In 1999, Alon observed that Chevalley's theorem follows from the Combinatorial Nullstellensatz he had just proved. On the other hand, this approach allowed Brink (2011) to prove a similar result in general fields $F$, but restricting the roots to belong to a product set $A_1\times\dots\times A_n$, where $A_1,\dots,A_n$ are finite subsets of $F$ of cardinality $q$. See that paper of Clarke, Forrow and Schmitt for further developments, in particular a version of Warning's second theorem.

b) In the case of hypersurfaces (with the notation of theorem 1, $m=1$), Ax proved in 1964 that the cardinality of $Z$ is divisible not only by $p$, but by $q$. This led to renewed interest in the following years, especially in the works of Katz, Esnault, Berthelot, and the well has not dried up yet.

c) In 2011, Heath-Brown published a paper where he uses Ax's result to strengthen the congruence modulo $p$ of proposition 4 to a congruence modulo $q$.

d) By a Weil restriction argument, a 1995 paper of Moreno-Moreno partially deduces the Chevalley-Warning theorem over a field of cardinality $q$ from its particular case over the prime field. I write partially because they obtain a divisibility by an expression of the form $p^{\lceil f \alpha\rceil}$, while one expects $q^{\lceil \alpha}=p^{f\lceil\alpha\rceil}$. However, the same argument allows them to obtain a stronger bound which does not involve not the degrees of the polynomials, but the $p$-weights of these degrees, that is the sum of their digits in their base $p$ expansions. Again, they obtain a divisibility by an expression of the form $p^{\lceil f\beta\rceil}$, and it is a natural question to wonder whether the divisibility by $p^{f\lceil\beta\rceil}$ can be proved.

Counting points and counting curves on varieties — Tribute to Daniel Perrin

$\require{enclose}\def\VarC{\mathrm{Var}_{\mathbf C}}\def\KVarC{K_0\VarC}$
Daniel Perrin is a French algebraic geometer who turned 70 last year. He his also well known in France for his wonderful teaching habilities. He was one of the cornerstones of the former École normale supérieure de jeunes filles, before it merged in 1985 with the rue d'Ulm school. From this time remains a Cours d'algèbre which is a must for all the students (and their teachers) who prepare the agrégation, the highest recruitment process for French high schools. He actually taught me Galois theory (at École normale supérieure in 1990/1991) and Algebraic Geometry (the year after, at Orsay). His teaching restlessly stresses  the importance of examples. He has also been deeply involved in training future primary school teachers, as well as in devising the mathematical curriculum of high school students: he was responsible of the report on geometry. It has been a great honor for me to be invited to lecture during the celebration of his achievements that took place at Orsay on November, 23, 2016.

Diophantine equations are a source of numerous arithmetic problems. One of them has been put forward by Manin in the 80s and consists in studying the behavior of the number of solutions of such equations of given size, when the bound grows to infinity. A geometric analogue of this question considers the space of all curves with given degree which are drawn on a fixed complex projective, and is interested in their behavior when the degree tends to infinity. This was the topic of my lecture and is the subject of this post.

Let us first begin with an old problem, apparently studied by Dirichlet around 1840, and given a rigorous solution by Chebyshev and Cesáro around 1880: the probability that two integers be coprime is equal to $6/\pi^2$. Of course, there is no probability on the integers that has the properties one would expect, such as being invariant by translation, and the classical formalization of this problem states that the numbers of pairs $(a,b)$ of integers such that $1\leq a,b\leq n$ and $\gcd(a,b)=1$ grows as $n^2 \cdot 6/\pi^2$ when $n\to+\infty$,

This can be proved relatively easily, for example as follows. Without the coprimality condition, there are $n^2$ such integers. Now one needs to remove those pairs both of which entries are multiples of $2$, and there are $\lfloor n/2\rfloor^2$ of those, those where $a,b$ are both multiples of $3$ ($\lfloor n/3\rfloor^2$), and then comes $5$, because we have already removed those even pairs, etc. for all prime numbers. But in this process, we have removed twice the pairs of integers both of which entries are multiples of $2\cdot 3=6$, so we have to add them back, and then remove the pairs of integers both of which are multiples of $2\cdot 3\cdot 5$, etc. This leads to the following formula for
the cardinality $C(n)$ we are interested in:

$\displaystyle
 C(n) = n^2 - \lfloor\frac n2\rfloor^2 - \lfloor \frac n3\rfloor^2-\lfloor \frac n5\rfloor^2 - \dots
+ \lfloor \frac n{2\cdot 3}\rfloor^2+\lfloor\frac n{2\cdot 5}\rfloor^2+\dots
- \lfloor \frac n{2\cdot 3\cdot 5} \rfloor^2 - \dots $.

Approximating $\lfloor n/a\rfloor$ by $n/a$, this becomes

$\displaystyle
C(n) \approx  n^2 - \left(\frac n2\right)-^2 - \left (\frac n3\right)^2-\left( \frac n5\rfloor\right)^2 - \dots
+ \left (\frac n{2\cdot 3}\right)^2+\left(\frac n{2\cdot 5}\right)^2+\dots
- \left (\frac n{2\cdot 3\cdot 5} \right)^2 - \dots $

which we recognize as

$\displaystyle
C(n)\approx n^2 \left(1-\frac1{2^2}\right) \left(1-\frac1{3^2}\right)\left(1-\frac1{5^2}\right) \dots
=n^2/\zeta(2)$,

where $\zeta(2)$ is the value at $s=2$ of Riemann's zeta function $\zeta(s)$. Now, Euler had revealed the truly arithmetic nature of $\pi$ by proving in 1734 that $\zeta(2)=\pi^2/6$. The approximations we made in this calculation can be justified, and this furnishes a proof of the above claim.

We can put this question about integers in a broader perspective if we recall that the ring $\mathbf Z$ is a principal ideal domain (PID) and study the analogue of our problem in other PIDs, in particular for $\mathbf F[T]$, where $\mathbf F$ is a finite field; set $q=\operatorname{Card}(\mathbf F)$. The above proof can be adapted easily (with simplifications, in fact) and shows that number of pairs $(A,B)$ of monic polynomials of degrees $\leq n$ such that $\gcd(A,B)=1$ grows as $q^n(1-1/q)$ when $n\to+\infty$. The analogy becomes stronger if one observes that $1/(1-1/q)$ is the value at $s=2$ of $1/(1-q^{1-s})$, the Hasse-Weil zeta function of the affine line over $\mathbf F$.

What can we say about our initial question if we replace the ring $\mathbf Z$ with the PID $\mathbf C[T]$? Of course, there's no point in counting the set of pairs $(A,B)$ of coprime monic polynomials of degree $\leq n$ in $\mathbf C[T]$, because this set is infinite. Can we, however, describe this set? For simplicity, we will consider here the set $V_n$ of pairs of coprime monic polynomials of degree precisely $n$. If we identify a monic polynomial of degree $n$ with the sequence of its coefficients, we then view $V_n$ as a subset of $\mathbf C^{n}\times\mathbf C^n$. We first observe that $V_n$ is an Zariski open subset of $\mathbf C^{2n}$: its complement $W_n$ is defined by the vanishing of a polynomial in $2n$ variables — the resultant of $A$ and $B$.

When $n=0$, we have $V_0=\mathbf C^0=\{\mathrm{pt}\}$.

Let's look at $n=1$: the polynomials $A=T+a$ and $B=T+b$ are coprime if and only if $a\neq b$;
consequently, $V_1$ is the complement of the diagonal in $\mathbf C^2$.

For $n=2$, this becomes more complicated: the resultant of the polynomials $T^2+aT+b$ and $T^2+cT+d$ is equal to $a^2d-abc-adc+b^2-2bd+bc^2+d^2$; however, it looks hard to guess some relevant properties of $V_n$ (or of its complement) just by staring at this equation. In any case, we can say that $V_2$ is the complement in $\mathbf C^4$ of the union of two sets, corresponding of the degree of the gcd of $(A,B)$. When $\gcd(A,B)=2$, one has $A=B$; this gives the diagonal, a subset of $\mathbf C^4$ isomorphic to $\mathbf C^2$; the set of pairs of polynomials $(A,B)$ whose gcd has degree $1$ is essentially $\mathbf C\times V_1$: multiply a pair $(A_1,B_1)$ of coprime polynomials of degree $1$ by an arbitrary polynomial of the form $(T-d)$.
Consequently,
\begin{align}V_2&=\mathbf C^4 - \left( \mathbf C^2 \cup \mathbf C\times V_1\right)\\
&= \mathbf C^4 - \left( \enclose{updiagonalstrike}{\mathbf C^2}\cup \left(\mathbf C\times (\mathbf C^2-\enclose{updiagonalstrike}{\mathbf C})\right)\right)\\
&=\mathbf C^4-\mathbf C^3
\end{align}
if we cancel the two $\mathbf C^2$ that appear. Except that this makes no sense!

However, there is a way to make this computation both meaningful and rigorous, and it consists in working in the Grothendieck ring $\KVarC$ of complex algebraic varieties. Its additive group is generated by isomorphism classes of algebraic varieties, with relations of the form $[X]=[U]+[Z]$ for every Zariski closed subset $Z$ of an algebraic variety $X$, with complement $U=X-Z$. This group has a natural ring structure for which $[X][Y]=[X\times Y]$. Its unit element is the class of the point, $[\mathbf A^0]$ if one wishes. An important element of this ring $\KVarC$ is the class $\mathbf L=[\mathbf A^1]$ of the affine line. The natural map $e\colon \VarC\to \KVarC$ given by $e(X)=[X]$ is the universal Euler characteristic: it is the universal map from $\VarC$ to a ring such that $e(X)=e(X-Z)+e(Z)$ and $e(X\times Y)=e(X)e(Y)$, where $X,Y$ are complex varieties and $Z$ is a Zariski closed subset of $X$.

In particular, it generalizes the classical Euler characteristic, the alternate sum of the dimensions of the cohomology groups (with compact support, if one wishes) of a variety. A subtler invariant of $\KVarC$ is given by mixed Hodge theory: there exists a unique ring morphism $\chi_{\mathrm H}\KVarC\to\mathbf Z[u,v]$ such that for every complex variety $X$, $\chi_{\mathrm H}([X])$ is the Hodge-Deligne polynomial of $X$. In particular, if $X$ is projective and smooth, $\chi_{\mathrm H}([X])=\sup_{p,q} \dim h^q(X,\Omega^p_X) u^pv^q$. If one replaces the field of complex numbers with a finite field $\mathbf F$, one may actually count the numbers of $\mathbf F$-points of $X$, and this furnishes yet another generalized Euler characteristic.

The preceding calculation shows that $e(V_0)=1$, $e(V_1)=\mathbf L^2-\mathbf L$ and $e(V_2)=\mathbf L^4-\mathbf L^3$; more generally, one proves by induction that $e(V_n)=\mathbf L^{2n}-\mathbf L^{2n-1}$ for every integer $n\geq 0$.

Equivalently, one has $e(W_n)=\mathbf L^{2n-1}$ for all $n$. I have to admit that I see no obvious reason for the class of $W_n$ to be equal to that of an affine space. However, as Ofer Gabber and Jean-Louis Colliot-Thélène pointed out to me during the talk, this resultant is the difference of two homogeneous polynomials $p-q$ of degrees $d=2$ and $d+1=3$; consequently, the locus it defines is a rational variety — given $a,b,c$, there is generically a unique $t$ such that $p-q$ vanishes at $(at,bt,ct,t)$.

These three results have a common interpretation if one brings in the projective line $\mathbf P_1$. Indeed, pairs $(a,b)$ of coprime integers (up to $\pm1$) correspond to rational points on $\mathbf P_1$, and if $\mathbf F$ is a field, then pairs $(A,B)$ of coprime polynomials in $\mathbf F[T]$ correspond (up to $\mathbf F^\times$) to elements of $\mathbf P_1(\mathbf F(T))$.
In both examples, the numerical datum $\max(|a|,|b|)$ or $\max(\deg(A),\deg(B))$ is called the height of the corresponding point.

In the case of the ring $\mathbf Z$, or in the case of the ring $\mathbf F[T]$ where $\mathbf F$ is a finite field, one has an obvious but fundamental finiteness theorem: there are only finitely many points of $\mathbf P_1$ with bounded height. In the latter case, $\mathbf C[T]$, this naïve finiteness does not hold. Nevertheless, if one sees $\mathbf P_1(\mathbf C(T))$ as an infinite dimensional variety — one needs infinitely many complex numbers to describe a rational function, then the points of bounded height constitute what is called a bounded family, a “finite dimensional” constructible set.

The last two examples have a common geometric interpretation. Namely, $\mathbf F(T)$ is the field of functions of a projective smooth algebraic curve $C$ over $\mathbf F$; in fact, $C$ is the projective line again, but we may better ignore this coincidence. Then a point $x\in\mathbf P_1(\mathbf F(T))$
corresponds to a morphism $\varepsilon_x\colon C\to\mathbf P_1$, and the formula $H(x)=\deg(\epsilon_x^*\mathscr O(1))$ relates the height $H(x)$ of $x$ to the degree of the morphism $\varepsilon_x$.

Since the notion of height generalizes from $\mathbf P_1$ to projective spaces $\mathbf P_n$ of higher dimension (and from $\mathbf Q$ to general number fields), this suggests a general question. Let $V\subset\mathbf P_n$ be a projective variety over a base field $k$ hat can one say about the set of points $x\in V(k)$ such that $H(x)\leq B$, when the bound $B$ grows to $\infty$?
The base field $k$ can be either a number field, or the field of functions $\mathbf F(C)$ of a curve $C$ over a finite field $\mathbf F$, or the field of functions $\mathbf C(C)$ of a curve over the complex numbers. In the last two cases, the variety can even be taken to be constant, deduced from a variety $V_0$ over $\mathbf F$ or $\mathbf C$.

1. When $k$ is a number field, this set is a finite set; how does its cardinality grows? This is a question that Batyrev and Manin have put forward at the end of the 80s, and which has attracted a lot of attention since.
2. When $k=\mathbf F(C)$ is a function field over a finite field, this set is again a finite set; how does its cardinality grows? This question has been proposed by Emmanuel Peyre by analogy with the question of Batyrev and Manin.
3. When $k=\mathbf C(C)$ is a function field over $\mathbf C$, this set identifies with a closed subscheme of the Grothendieck-Hilbert scheme of $V$; what can one say about its geometry, in particular about its class in $\KVarC$? Again, this question has been proposed by Emmanuel Peyre around 2000.

In a forthcoming post, I shall recall some results on these questions, especially the first one, and in particular explain an approach based on the Fourier summation formula. I will then explain a theorem proved with François Loeser where we make use of Hrushovski–Kazhdan's motivic Fourier summation formula in motivic integration to prove an instance of the third question.

Weierstrass's approximation theorem

I had to mentor an Agrégation leçon entitled Examples of dense subsets. For my own edification (and that of the masses), I want to try to record here as many proofs as of the Weierstrass density theorem as I can : Every complex-valued continuous function on the closed interval $[-1;1]$ can be uniformly approximated by polynomials. I'll also include as a bonus the trigonometric variant: Every complex-valued continuous and $2\pi$-periodic function on $\mathbf R$ can be uniformly approximated by trigonometric polynomials.

1. Using the Stone theorem.

This 1937—1948 theorem is probably the final conceptual brick to the edifice of which Weierstrass laid the first stone in 1885. It asserts that a subalgebra of continuous functions on a compact totally regular (e.g., metric) space is dense for the uniform norm if and only if it separates points. In all presentations that I know of, its proof requires to establish that the absolute value function can be uniformly approximated by polynomials on $[-1;1]$:

• Stone truncates the power series expansion of the function \[ x\mapsto \sqrt{1-(1-x^2)}=\sum_{n=0}^\infty \binom{1/2}n (x^2-1)^n, \] bounding by hand the error term.
• Bourbaki (Topologie générale, X, p. 36, lemme 2) follows a more elementary approach and begins by proving  that the function $x\mapsto \sqrt x$ can be uniformly approximated by polynomials on $[0;1]$. (The absolute value function is recovered since $\mathopen|x\mathclose|\sqrt{x^2}$.) To this aim, he introduces the sequence of polynomials given by $p_0=0$ and $p_{n+1}(x)=p_n(x)+\frac12\left(x-p_n(x)^2\right)$ and proves by induction the inequalities \[ 0\leq \sqrt x-p_n(x) \leq \frac{2\sqrt x}{2+n\sqrt x} \leq \frac 2n\] for $x\in[0;1]$ and $n\geq 0$. This implies the desired result.

The algebra of polynomials separates points on the compact set $[-1;1]$, hence is dense. To treat the case of trigonometric polynomials, consider Laurent polynomials on the unit circle.

2. Convolution.

Consider an approximation $(\rho_n)$ of the Dirac distribution, i.e., a sequence of continuous, nonnegative and compactly supported functions on $\mathbf R$ such that $\int\rho_n=1$ and such that for every $\delta>0$, $\int_{\mathopen| x\mathclose|>\delta} \rho_n(x)\,dx\to 0$. Given a continuous function $f$ on $\mathbf R$, form the convolutions defined by $f*\rho_n(x)=\int_{\mathbf R} \rho_n(t) f(x-t)\, dt$. It is classical that $f*\rho_n$ converges uniformly on every compact to $f$.

Now, given a continuous function $f$ on $[-1;1]$, one can extend it to a continuous function with compact support on $\mathbf R$ (defining $f$ to be affine linear on $[-2;-1]$ and on $[1;2]$, and to be zero outside of $[-2;2]$. We want to choose $\rho_n$ so that $f*\rho_n$ is a polynomial on $[-1;1]$. The basic idea is just to choose a parameter $a>0$, and to take $\rho_n(x)= c_n (1-(x/a)^2)^n$ for $\mathopen|x\mathclose|\leq a$ and $\rho_n(x)=0$ otherwise, with $c_n$ adjusted so that $\int\rho_n=1$. Let us write $f*\rho_n(x)=\int_{-2}^2 \rho_n(x-t) f(t)\, dt$; if $x\in[-1;1]$ and $t\in[-2:2]$, then $x-t\in [-3;3]$ so we just need to be sure that $\rho_n$ is a polynomial on that interval, which we get by taking, say, $a=3$. This shows that the restriction of $f*\rho_n$ to $[-1;1]$ is a polynomial function, and we're done.

This approach is more or less that of D. Jackson (“A Proof of Weierstrass's Theorem,” Amer. Math. Monthly, 1934). The difference is that he considers continuous functions on a closed interval contained in $\mathopen]0;1\mathclose[$ which he extends linearly to $[0;1]$ so that they vanish at $0$ and $1$; he considers the same convolution, taking the parameter $a=1$.

Weierstrass's own proof (“Über die analytische Darstellbarkeit sogenannter willkurlicher Functionen einer reellen Veranderlichen Sitzungsberichteder,” Königlich Preussischen Akademie der Wissenschaften zu Berlin, 1885) was slightly more sophisticated: he first showed approximation by convolution with the Gaussian kernel  defined by $ \rho_n(t) =\sqrt{ n} e^{- \pi n t^2}$, and then expanded the kernel as a power series, a suitable truncation of which furnishes the desired polynomials.

As shown by Jacskon, the same approach works easily (in a sense, more easily) for $2\pi$-periodic functions, considering the kernel defined by $\rho_n(x)=c_n(1+\cos(x))^n$, where $c_n$ is chosen so that \int_{-\pi}^\pi \rho_n=1$.

3. Bernstein polynomials.

Take a continuous function $f$ on $[0;1]$ and, for $n\geq 0$, set \[ B_nf(x) = \sum_{k=0}^n f(k/n) \binom nk t^k (1-t)^{n-k}.\] It is classical that $B_nf$ converges uniformly to $f$ on $[0;1]$.

There are two classical proofs of Bernstein's theorem. One is probabilistic and consists in observing that $B_nf(x)$ is the expected value of $f(S_n)$, where $S_n$ is the sum of $n$ i.i.d. Bernoulli random variables with parameter $x\in[0;1]$. Another (generalized as the Korovkin theorem, “On convergence of linear positive operators in the space of continuous functions”, Dokl. Akad. Nauk SSSR (N.S.), vol. 90,‎ 1953) consists in showing (i) that for $f=1,x,x^2$, $B_nf$ converges uniformly to $f$ (an explicit calculation), (ii) that if $f\geq 0$, then $B_nf\geq 0$ as well, (iii) for every $x\in[0;1]$, squeezing $f$ inbetween two quadratic polynomials $f^+$ and $f_-$ such that $f^+(x)-f^-(x)$ is as small as desired.

A trigonometric variant would be given by Fejér's theorem that the Cesàro averages of a Fourier series of a continuous, $2\pi$-periodic function converge uniformly to that function. In turn, Fejér's theorem can be proved in both ways, either by convolution (the Fejér kernel is nonnegative), or by a Korovkine-type argument (replacing $1,x,x^2$ on $[0;1]$ by $1,z,z^2,z^{-1},z^{-2}$ on the unit circle).


4. Using approximation by step functions.

This proof originates with a paper of H. Kuhn, “Ein elementarer Beweis des Weierstrasschen Approximationsatzes,” Arch. Math. 15 (1964), p. 316–317.

Let us show that for every $\delta\in\mathopen]0,1\mathclose[$ and every $\varepsilon>0$, there exists a polynomial $p$ satisfying the following properties:

• $0\leq p(x)\leq \varepsilon$ for $-1\leq x\leq-\delta$;
• $0\leq p(x)\leq 1$ for $-\delta\leq x\leq \delta$;
• $1-\varepsilon\leq p(x)\leq 1$ for $\delta\leq x\leq 1$.

In other words, these polynomials approximate the (discontinuous) function $f$ on $[-1;1]$ defined by $f(x)=0$ for $x< 0$, $f(x)=1$ for $x> 0$ and $f(0)=1/2$.

A possible formula is $p(x)=(1- ((1-x)/2))^n)^{2^n}$, where $n$ is a large enough integer. First of all, one has $0\leq (1-x)/2\leq 1$ for every $x\in[-1;1]$, so that $0\leq p(x)\leq 1$. Let $x\in[-1;-\delta]$; then one has $(1-x)/2\geq (1+\delta)/2$, hence $p(x)\leq (1-((1+\delta)/2)^n)^{2^n}$, which can be made arbitrarily small when $n\to\infty$. Let finally $x\in[\delta;1]$; then $(1-x)/2\geq (1-\delta)/2$, hence $p(x)\geq (1-((1-\delta)/2)^n)^{2^n}\geq 1- (1-\delta)^n$, which can be made arbitrarily close to $1$ when $n\to\infty$.

By translation and dilations, the discontinuity can be placed at any element of $[0;1]$. Let now $f$ be an arbitrary step function and let us write it as a linear combination $f=\sum a_i f_i$, where $f_i$ is a $\{0,1\}$-valued step function. For every $i$, let $p_i$ be a polynomial that approximates $f_i$ as given above. The linear combination $\sum a_i p_i$ approximates $f$ with maximal error $\sup(\mathopen|a_i\mathclose|)$.

Using uniform continuity of continuous functions on $[-1;1]$, every continuous function can be uniformly approximated by a step function. This concludes the proof.

5. Using approximation by piecewise linear functions.

As in the proof of Stone's theorem, one uses the fact that the function $x\mapsto \mathopen|x\mathclose|$ is uniformly approximated by a sequence of polynomial on $[-1;1]$. Consequently,  so are the functions $x\mapsto \max(0,x)=(x+\mathopen|x\mathclose|)/2 $ and $x\mapsto\min(0,x)=(x-\mathopen|x\mathclose|)/2$. By translation and dilation, every continuous piecewise linear function on $[-1;1]$ with only one break point is uniformly approximated by polynomials. By linear combination, every continuous piecewise linear affine function is uniformly approximated by polynomials.
By uniform continuity, every continuous function can be uniformly approximated by continuous piecewise linear affine functions. Weierstrass's theorem follows.

6. Moments.

A linear subspace $A$ of a Banach space is dense if and only if every continuous linear form which vanishes on $A$ is identically $0$. In the present case, the dual of $C^0([-1;1],\mathbf C)$ is the space of complex measures on $[-1;1]$ (Riesz theorem, if one wish, or the definition of a measure). So let $\mu$ be a complex measure on $[-1;1]$ such that $\int_{-1}^1 t^n \,d\mu(t)=0$ for every integer $n\geq 0$; let us show that $\mu=0$. This is the classical problem of showing that a complex measure on $[-1;1]$ is determined by its moments. In fact, the classical proof of this fact runs the other way round, and there must exist ways to reverse the arguments.

One such solution is given in Rudin's Real and complex analysis, where it is more convenient to consider functions on the interval $[0;1]$. So, let $F(z)=\int_0^1 t^z \,d\mu(t)$. The function $F$ is holomorphic and bounded on the half-plane $\Re(z)> 0$ and vanishes at the positive integers. At this point, Rudin makes a conform transformation to the unit disk (setting $w=(z-1)/(z+1)$) and gets a  bounded function on the unit disk with zeroes at $(n-1)/(n+1)=1-2/(n+1)$, for $n\in\mathbf N$, and this contradicts the fact that the series $\sum 1/(n+1)$ diverges.

In Rudin, this method is used to prove the more general Müntz–Szász theorem according to which the family $(t^{\lambda_n})$ generates a dense subset of $C([0;1])$ if and only if $\sum 1/\lambda_n=+\infty$.

Here is another solution I learnt in a paper by L. Carleson (“Mergelyan's theorem on uniform polynomial approximation”, Math. Scand., 1964).

For every complex number $a$ such that $\mathopen|a\mathclose|>1$, one can write $1/(t-a)$ as a converging power series. By summation, this quickly gives that
\[ F(a) = \int_{-1}^1 \frac{1}{t-a}\, d\mu(t) \equiv 0. \]
Observe that this formula defines a holomorphic function on $\mathbf C\setminus[-1;1]$; by analytic continuous, one thus has $F(a)=0$ for every $a\not\in[-1;1]$.
Take a $C^2$-function $g$ with compact support on the complex plane. For every $t\in\mathbf C$, one has the following formula
\[ \iint \bar\partial g(z) \frac{1}{t-z} \, dx\,dy = g(t), \]
which implies, by integration and Fubini, that
\[ \int_{-1}^1 g(t)\,d\mu(t) = \iint \int \bar\partial g(z) \frac1{t-z}\,d\mu(t)\,dx\,dy = \iint \bar\partial g(z) F(z)\,dx\, dy= 0. \]
On the other hand, every $C^2$ function on $[-1;1]$ can be extended to such a function $g$, so that the measure $\mu$ vanishes on every $C^2$ function on $[-1;1]$. Approximating a continuous function by a $C^2$ function (first take a piecewise linear approximation, and round the corners), we get that $\mu$ vanishes on every continuous function, as was to be proved.

7. Chebyshev/Markov systems.

This proof is due to P. Borwein and taken from the book Polynomials and polynomial inequalities, by P. Borwein and T. Erdélyi (Graduate Texts in Maths, vol. 161, 1995). Let us say that a sequence $(f_n)$ of continuous functions on an interval $I$ is a Markov system (resp. a weak Markov system) if for every integer $n$, every linear combination of $(f_0,\dots,f_n)$ has at most $n$ zeroes (resp. $n$ sign changes) in $I$.

Given a Markov system $(f_n)$, one defines a sequence $(T_n)$, where $T_n-f_n$ is the element of $\langle f_0,\dots,f_{n-1}\rangle$ which is the closest to $f_n$. The function $T_n$ has $n$ zeroes on the interval $I$; let $M_n$ be the maximum distance between two consecutive zeroes.

Borwein's theorem  (Theorem 4.1.1 in the mentioned book) then asserts that if the sequence $(f_n)$ is a Markov system consisting of $C^1$ functions, then its linear span is dense in $C(I)$ if and only if $M_n\to 0$.

The sequence of monomials $(x^n)$ on $I=[-1;1]$ is of course a Markov system.  In this case, the polynomial $T_n$ is the $n$th Chebyshev polynomial, given by $T_n(2\cos(x))=2\cos(nx)$, and its roots are given by $2\cos((\pi+2k\pi)/2n)$, for $k=0,\dots,n-1$, and $M_n\leq \pi/n$. This gives yet another proof of Weierstrass's approximation theorem.